Stable Diffusion v3.5 Large
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A state-of-the-art text-to-image generative model designed to create high-resolution images based on textual prompts. It excels in producing diverse and high-quality outputs, making it suitable for professional applications.
If you don’t have an API key for the AI/ML API yet, feel free to use our .
Let's generate an image using a simple prompt.
We obtained the following 1024x576 image by running this code example:
"A T-Rex relaxing on a beach, lying on a sun lounger and wearing sunglasses."
"A highly detailed T-Rex relaxing on a sunny beach, lying on a wooden sun lounger and wearing stylish sunglasses. Its skin is covered in realistic, finely textured scales with natural color variations — rough and weathered like that of large reptiles. Sunlight reflects subtly off the individual scales. The background includes palm trees, gentle waves, and soft sand partially covering the T-Rex's feet. The scene is rendered with cinematic lighting and a natural color palette."
"num_inference_steps": 40
"Racoon eating ice-cream"
"A T-Rex relaxing on a beach, lying on a sun lounger and wearing sunglasses. Vector illustration style. Top-down view, with visible palm trees, seagulls, and a strip of water."
"num_inference_steps": 40
square_hd
The size of the generated image.
The description of elements to avoid in the generated image.
The CFG (Classifier Free Guidance) scale is a measure of how close you want the model to stick to your prompt when looking for a related image to show you.
The number of inference steps to perform.
If set to True, the safety checker will be enabled.
true
The text prompt describing the content, style, or composition of the image to be generated.
The number of images to generate.
1
The same seed and the same prompt given to the same version of the model will output the same image every time.
The format of the generated image.
jpeg
Possible values: No content